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Improved techniques for training Score-Based generative models

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13 Pith papers citing it
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abstract

Score-based generative models can produce high quality image samples comparable to GANs, without requiring adversarial optimization. However, existing training procedures are limited to images of low resolution (typically below 32x32), and can be unstable under some settings. We provide a new theoretical analysis of learning and sampling from score models in high dimensional spaces, explaining existing failure modes and motivating new solutions that generalize across datasets. To enhance stability, we also propose to maintain an exponential moving average of model weights. With these improvements, we can effortlessly scale score-based generative models to images with unprecedented resolutions ranging from 64x64 to 256x256. Our score-based models can generate high-fidelity samples that rival best-in-class GANs on various image datasets, including CelebA, FFHQ, and multiple LSUN categories.

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representative citing papers

Denoising Diffusion Implicit Models

cs.LG · 2020-10-06 · unverdicted · novelty 8.0

DDIMs construct non-Markovian diffusion processes that share DDPM training objectives but allow much faster reverse sampling, demonstrated empirically at 10-50x wall-clock speedup.

DiffWave: A Versatile Diffusion Model for Audio Synthesis

eess.AS · 2020-09-21 · unverdicted · novelty 8.0

DiffWave is a non-autoregressive diffusion model that generates high-fidelity audio waveforms from noise in constant steps, matching WaveNet vocoder quality while being orders of magnitude faster and outperforming prior models in unconditional generation.

Denoising Diffusion Probabilistic Models

cs.LG · 2020-06-19 · accept · novelty 8.0

Denoising diffusion probabilistic models generate high-quality images by learning to reverse a fixed forward diffusion process, achieving FID 3.17 on CIFAR10.

Diffusion Models Beat GANs on Image Synthesis

cs.LG · 2021-05-11 · accept · novelty 7.0

Diffusion models with architecture improvements and classifier guidance achieve superior FID scores to GANs on unconditional and conditional ImageNet image synthesis.

Improved Denoising Diffusion Probabilistic Models

cs.LG · 2021-02-18 · accept · novelty 7.0

Targeted tweaks to DDPMs produce competitive likelihoods and high-quality samples, with learned reverse variances enabling 10x faster sampling and predictable scaling with compute.

Stable Video Diffusion: Scaling Latent Video Diffusion Models to Large Datasets

cs.CV · 2023-11-25 · conditional · novelty 6.0

Stable Video Diffusion scales latent video diffusion models via text-to-image pretraining, video pretraining on curated data, and high-quality finetuning to produce competitive text-to-video and image-to-video results while enabling motion LoRA and multi-view 3D applications.

Exploring the flavor structure of leptons via diffusion models

hep-ph · 2025-03-27 · unverdicted · novelty 5.0

Applies diffusion models to generate 10,000 neutrino mass matrices consistent with oscillation parameters in a seesaw model, revealing non-trivial distributions in CP phases and 0νββ effective mass.

citing papers explorer

Showing 13 of 13 citing papers.

  • Denoising Diffusion Implicit Models cs.LG · 2020-10-06 · unverdicted · none · ref 21

    DDIMs construct non-Markovian diffusion processes that share DDPM training objectives but allow much faster reverse sampling, demonstrated empirically at 10-50x wall-clock speedup.

  • DiffWave: A Versatile Diffusion Model for Audio Synthesis eess.AS · 2020-09-21 · unverdicted · none · ref 13

    DiffWave is a non-autoregressive diffusion model that generates high-fidelity audio waveforms from noise in constant steps, matching WaveNet vocoder quality while being orders of magnitude faster and outperforming prior models in unconditional generation.

  • Denoising Diffusion Probabilistic Models cs.LG · 2020-06-19 · accept · none · ref 56

    Denoising diffusion probabilistic models generate high-quality images by learning to reverse a fixed forward diffusion process, achieving FID 3.17 on CIFAR10.

  • Hierarchical Text-Conditional Image Generation with CLIP Latents cs.CV · 2022-04-13 · accept · none · ref 49

    A hierarchical prior-decoder model using CLIP latents generates more diverse text-conditional images than direct methods while preserving photorealism and caption fidelity.

  • SDEdit: Guided Image Synthesis and Editing with Stochastic Differential Equations cs.CV · 2021-08-02 · conditional · none · ref 11

    SDEdit performs guided image synthesis and editing by adding noise to inputs and refining them via denoising with a diffusion model's SDE prior, outperforming GAN methods in human studies without task-specific training.

  • Diffusion Models Beat GANs on Image Synthesis cs.LG · 2021-05-11 · accept · none · ref 58

    Diffusion models with architecture improvements and classifier guidance achieve superior FID scores to GANs on unconditional and conditional ImageNet image synthesis.

  • Improved Denoising Diffusion Probabilistic Models cs.LG · 2021-02-18 · accept · none · ref 8

    Targeted tweaks to DDPMs produce competitive likelihoods and high-quality samples, with learned reverse variances enabling 10x faster sampling and predictable scaling with compute.

  • Diffusion Models for Sampling Near Criticality in Lattice Field Theories hep-lat · 2026-07-09 · accept · none · ref 59 · internal anchor

    Fully convolutional diffusion models trained on small lattices transfer to unseen larger volumes for 2D/3D phi^4 sampling across phases, matching or beating same-size training on most observables.

  • GenSBI: Generative Methods for Simulation-Based Inference in JAX cs.LG · 2026-05-26 · unverdicted · none · ref 59

    GenSBI delivers JAX-native implementations of generative SBI methods with transformer backbones and reports near-ideal calibration scores on standard benchmarks.

  • Stable Video Diffusion: Scaling Latent Video Diffusion Models to Large Datasets cs.CV · 2023-11-25 · conditional · none · ref 86

    Stable Video Diffusion scales latent video diffusion models via text-to-image pretraining, video pretraining on curated data, and high-quality finetuning to produce competitive text-to-video and image-to-video results while enabling motion LoRA and multi-view 3D applications.

  • Knowledge Distillation in Iterative Generative Models for Improved Sampling Speed cs.LG · 2021-01-07 · unverdicted · none · ref 36

    Denoising Student distills the multi-step denoising process of score-based and diffusion models into a single forward pass, matching GAN sampling speed while producing comparable sample quality on CIFAR-10, CelebA, and 256x256 LSUN.

  • Exploring the flavor structure of leptons via diffusion models hep-ph · 2025-03-27 · unverdicted · none · ref 40

    Applies diffusion models to generate 10,000 neutrino mass matrices consistent with oscillation parameters in a seesaw model, revealing non-trivial distributions in CP phases and 0νββ effective mass.

  • A Unified Measure-Theoretic View of Diffusion, Score-Based, and Flow Matching Generative Models cs.LG · 2026-05-07 · unverdicted · none · ref 13

    Diffusion, score-based, and flow matching models are unified as instances of learning time-dependent vector fields inducing marginal distributions governed by continuity and Fokker-Planck equations.